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models
Stability AI Models
Stability AI is the lab behind Stable Diffusion, the open-weight model family that made high-quality image generation broadly accessible, and it is the second-largest partner catalogue on Segmind. The collection spans the full lineage: Stable Diffusion 3.5 and SD3 for current-generation text-to-image quality, SDXL and SDXL Turbo for the widely-adopted 1024px workhorse and its distilled fast variant, and the SD 1.5 and 2.1 pipelines that remain the base for an enormous ecosystem of community LoRAs and fine-tunes. Alongside the base models sit the control and editing pipelines that make Stable Diffusion practical for production work: ControlNet variants for pose, depth, canny-edge and scribble guidance, inpainting and outpainting models for targeted region edits, image-to-image for style-preserving transformation, and Stable Video Diffusion for animating stills. Stability's open-weight approach is what makes this collection distinctive: the models are transparent, extensively documented, and backed by the largest fine-tuning community in generative imaging, so a technique you find in a paper or a community post is usually reproducible here. On Segmind, every Stability model is a pay-per-use API endpoint with no GPU to provision. Chain them with upscalers, face restoration and background removal in Segmind Workflows to build complete automated image production pipelines.
Stable Diffusion 3.5 Turbo Text to Image
Stable Diffusion 3.5 Turbo offers exceptional customizability, efficient performance on consumer hardware, and diverse image outputs that accurately represent different skin tones and features, all while maintaining high-quality results and strong prompt adherence.
Stable Diffusion 3.5 Large Text to Image
Stable Diffusion 3.5 Large offers exceptional customizability, efficient performance on consumer hardware, and diverse image outputs that accurately represent different skin tones and features, all while maintaining high-quality results and strong prompt adherence.
Realdream Pony V9
Real Dream Pony V9 is an advanced image generation model based on the Stable Diffusion XL (SDXL) architecture, excelling in photorealism.
RealDream Lightning
RealDream is a sophisticated image generation model utilizing SDXL Lightning architecture. It creates incredibly realistic images from textual prompts. With the ability to excellently generate human portraits from the user's descriptive text.
Stable Diffusion 3 Medium Image to Image
Stable Diffusion 3 Medium image-to-image is a cutting-edge AI tool that uses advanced image-to-image technology to transform one image into another.
SD3 Medium Tile Controlnet
SD3 Medium Tile ControlNet is a large generative image model designed for generating detailed images based on textual prompts and tile-based input images.
SD3 Medium Canny Controlnet
Stable Diffusion 3 (SD3) Medium Canny ControlNet uses Canny edge detection to provide fine-grained control over the generated outputs.
SD3 Medium Pose Controlnet
Stable Diffusion 3 (SD3) Pose ControlNet is a large generative image model tailored for generating images based on text prompts while using pose information as guidance.
Motion Control SVD
Motion Control SVD is an innovative deep learning framework that breathes life into static images. By intelligently managing both camera and object motion, it empowers creators to achieve precise animation effects.
SDXL Img2Img
SDXL Img2Img is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers
SDXL Controlnet
SDXL ControlNet gives unprecedented control over text-to-image generation. SDXL ControlNet models Introduces the concept of conditioning inputs, which provide additional information to guide the image generation process
Stable Diffusion 3 Medium Text to Image
Stable Diffusion is a type of latent diffusion model that can generate images from text. It was created by a team of researchers and engineers from CompVis, Stability AI, and LAION. Stable Diffusion v2 is a specific version of the model architecture. It utilizes a downsampling-factor 8 autoencoder with an 865M UNet and OpenCLIP ViT-H/14 text encoder for the diffusion model. When using the SD 2-v model, it produces 768x768 px images. It uses the penultimate text embeddings from a CLIP ViT-H/14 text encoder to condition the generation process.
NewReality Lightning SDXL
NewReality Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
DreamShaper Lightning SDXL
DreamShaper Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Colossus Lightning SDXL
Colossus Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Samaritan Lightning SDXL
Samaritan Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Realism Lightning SDXL
Realism Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
ProtoVision Lightning SDXL
ProtoVision Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
NightVis Lightning SDXL
NightVis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
WildCard Lightning SDXL
WildCard Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Dynavis Lightning SDXL
Dynavis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Juggernaut Lightning SDXL
Juggernaut Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Realvis Lightning SDXL
Realvis Lightning SDXL is a lightning-fast text-to-image generation model. It can generate high-quality 1024px images in a few steps.
Samaritan 3D XL
Samaritan 3D XL leverages the robust capabilities of the SDXL framework, ensuring high-quality, detailed 3D character renderings.
Stable Video Diffusion
Takes image as input and returns a video.
IP-adapter Depth XL
IP Adapter Depth XL is built on the SDXL framework. This model integrates the IP Adapter and Depth preprocessor to offer unparalleled control and guidance in creating context-rich images.
SDXL Inpaint
This model is capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask
SSD Img2Img
This model uses SSD-1B to generate images by passing a text prompt and an initial image to condition the generation
SDXL-Openpose
This model leverages SDXL to generate the images with ControlNet conditioned on Human Pose Estimation.
SSD-Depth
This model leverages SSD-1B to generate the images with ControlNet conditioned on Depth Estimation
SSD-1B
SSD-1B efficiently generates high-quality, diverse images from text prompts in real-time.
Copax Timeless SDXL
The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.
Zavychroma SDXL
The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.
Realvis SDXL
The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.
Dreamshaper SDXL
The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software.
Word2img
Create beautifully designed words using Segmind’s word to image for your marketing purposes
Stable Diffusion Inpainting
Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask
Stable Diffusion img2img
This model uses diffusion-denoising mechanism as first proposed by SDEdit, Stable Diffusion is used for text-guided image-to-image translation. This model uses the weights from Stable Diffusion to generate new images from an input image using StableDiffusionImg2ImgPipeline from diffusers
Stable Diffusion XL 1.0
The SDXL model is the official upgrade to the v1.5 model. The model is released as open-source software
Reliberate
This model corresponds to the Stable Diffusion Reliberate checkpoint for detailed images at the cost of a super detailed prompt
Realistic Vision
This model corresponds to the Stable Diffusion Realistic Vision checkpoint for detailed images at the cost of a super detailed prompt
SD Outpainting
Stable Diffusion Outpainting can extend any image in any direction
Juggernaut Final
The most versatile photorealistic model that blends various models to achieve the amazing realistic images.
Epic Realism
This model corresponds to the Stable Diffusion Epic Realism checkpoint for detailed images at the cost of a super detailed prompt
Edge of Realism
This model corresponds to the Stable Diffusion Edge of Realism checkpoint for detailed images at the cost of a super detailed prompt
Cyber Realistic
The most versatile photorealistic model that blends various models to achieve the amazing realistic images.
ControlNet Soft Edge
This model corresponds to the ControlNet conditioned on Soft Edge.
ControlNet Scribble
This model corresponds to the ControlNet conditioned on Scribble images.
ControlNet Depth
This model corresponds to the ControlNet conditioned on Depth estimation.
ControlNet Canny
This model corresponds to the ControlNet conditioned on Canny edges.